
Explore how to generate high-quality images with Stable Diffusion, from basics and prompt engineering to fine-tuning, image-to-image, inpainting, and ControlNet, using Google Colab and open-source tools.
Explore stable diffusion, a latent diffusion-based deep learning model for text-guided image generation that creates high-quality images from prompts and supports image-to-image, inpainting, fine-tuning, super-resolution, and outpainting.
Explore stable diffusion fundamentals, including unconditional image generation, text-to-image prompts, and image-to-image conditioning, plus advantages like open source access and fast latent-space generation.
Explore the intuition of stable diffusion: transforming random noise into images via latent diffusion model and autoencoders. Understand the roles of U-NET, text-encoder, and CLIP in turning prompts into visuals.
Explore how stable diffusion combines gaussian noise latent seeds, a user prompt, CLIP text embeddings, and U-Net denoising, then decodes with a variational autoencoder to produce an image.
Learn to test stable diffusion via websites and Gradio interfaces, perform text-to-image and image-to-image tasks with inpainting, and tune sampling methods, steps, cfg scale, and seeds under Colab plan limitations.
Install and configure the stable diffusion pipeline in Google Colab by installing diffusers 0.11.1, accelerate, transformers, and memory-saving tools, set GPU and float16, and review commercial licenses.
Describe the subject in detail and use powerful keywords to define the final style, since prompts are text instructions that condition Stable Diffusion to generate a defined image.
Generate the first image by creating a simple prompt and running the pipeline. View the result at position zero and explore prompts, like a photograph of an apple.
Generate three images from a prompt using a grid_img function to arrange them in a PIL RGB grid with a configurable scale.
Explore how the seed initializes random numbers to create initial latent noise in stable diffusion and guarantee reproducible images when using the same seed.
Learn how inference steps, or denoising steps, govern image quality and generation speed in Stable Diffusion by testing values from 10 to 50 and noting scheduler effects.
Tune the CFG guidance scale to control how much the prompt shapes diffusion conditioning and image quality, with a default 7.5 and a 5–9 range.
Learn how negative prompts guide stable diffusion to exclude unwanted features, like mustaches, by using separate negative embeddings and refining prompts for higher quality images.
Learn how to implement negative prompts in Stable Diffusion to control outputs; generate three images of an old car, compare grayscale vs color, and visualize results in a grid.
Understand how Stable Diffusion models use pre-trained weights and training data to generate images, compare base models like v1.4–v2.1, and explain text embeddings, negative prompts, and resolutions.
Discover stable diffusion xl, with improved realism, face generation, and text-in-image capabilities, and learn about fast models like turbo and lightning, including how to compare models for best results.
Experiment with Stable Diffusion versions 1.5 and 2.x by loading pretrained pipelines, generating images with prompts and negative prompts, and evaluating quality differences.
Explore how fine-tuned models unlock specific styles in Stable Diffusion, load prebuilt models with a pipeline, and generate modern Disney style images using targeted prompts.
Explore how to switch the scheduler in a stable diffusion pipeline, compare PNDMScheduler and DDIMScheduler, and observe how different algorithms affect a terrier on the beach image.
Learn how training data shapes stable diffusion image generation and biases in occupations, skin tone, and gender, with Bloomberg and Washington Post insights, and mitigation through diverse data.
Prepare the environment for stable diffusion by configuring a Google Colab notebook with GPU acceleration, installing diffusers and PyTorch, and building the pipeline for repeatable results.
Craft precise prompts for stable diffusion by defining the subject, action, location, and image type; use a seed and a generator to produce varied outputs.
Apply modernist, impressionist, and realistic styles with color keywords in Stable Diffusion, and see how artist names like Van Gogh, Monet, and Botticelli shape outputs.
Explore how resolution shapes sharpness and detail with keywords like Unreal Engine and V-Ray. Tune prompts with cinematic lighting, ring lighting, and sunset, plus assets from Unsplash, Pixabay, and Pixiv.
Learn to use negative prompts in stable diffusion to remove undesired elements, fix anatomy, sharpen quality, and tailor outputs with parameters like negative_prompt, seeds, and testing different prompts.
Explore Stable Diffusion v2 image generation, compare 2.1 with earlier versions, apply the Euler Ancestral Discrete Scheduler on GPU, and leverage negative prompts to improve quality.
Explore generating arts and photographs with stable diffusion by testing prompts, steps, schedulers, and negative prompts to shape oil painting portraits and realistic images.
Generate landscapes and 3d images using prompts and a pipeline, adjusting height, width, steps, and negative prompts to produce high quality, realistic renders with blender and vray.
Generate drawings and architectural visuals with stable diffusion by crafting prompts and negative prompts, producing 8k sticker art and cinematic cityscapes like rio de janeiro while avoiding watermarks.
Test custom, fine-tuned models for stable diffusion and compare styles from Anything to Mitsua Diffusion One, while learning to manage safety with negative prompts.
Learn fine-tuning of text-to-image models with dreambooth and embedding, injecting a custom subject into Stable Diffusion using only three to five images and diversified datasets.
Set up a Google Colab workspace to fine-tune a custom model with your own images, including GPU, saving a copy in Drive, installing the diffusers library, and preparing for training.
Train a custom stable diffusion model with dreambooth by configuring instance and class prompts, organizing data folders, creating concepts_list.json, and running a ten-image training with set parameters.
Prepare the Stable Diffusion Dreambooth training environment by accessing weight folders, visualizing sample images, and converting trained weights to a ckpt file for image generation.
Learn to generate images from a trained Stable Diffusion model using a ckpt file and StableDiffusionPipeline on GPU, refining prompts, seeds, and saving results to disk.
Master AI image generation with stable diffusion by diagnosing common issues in fine-tuned models, optimizing prompts, training steps, image diversity, and inference settings for higher quality results.
Prepare the environment for image-to-image generation by installing libraries, loading the pretrained stable diffusion 1.5 pipeline on CUDA, and using a 512 by 512 Apple.jpg as guidance.
Create a prompt and run the pipeline with a manual seed to generate a photograph of an apple from an initial image; vary seeds for different results.
Explore the strength parameter in stable diffusion, which controls noise and variation from the input image; higher values produce more real-world outputs, lower values stay closer to the input.
Experiment with multiple image styles to transform an apple scene, applying watercolor and oil painting styles and varying prompts to preserve and alter the initial image.
Explore image-to-image generation with other models like modi in modern Disney style, experimenting with prompts, initial images, seeds, strength, and guidance_scale to transform landscapes into castles and cartoons.
Explore image-to-image editing with instruct pix2pix using stable diffusion pipelines to modify images with prompts, from adding mustaches to turning statues into cyborgs or in the style of van Gogh.
Learn to implement inpainting with stable diffusion in Google Colab, activate gpu acceleration, optimize memory, and remove objects using a mask and two images with a visualization grid.
Learn to use inpainting with Stable Diffusion by supplying an original image and a mask, setting a seed for repeatable results, and generating altered images like removing a dog.
Experiment with inpainting using Stable Diffusion by replacing a bench scene with bench02_img and bench02_img_mask in Google Colab, visualizing with grid_img and adjusting width, height, and prompts.
Explore controlnet to guide image generation with depth-to-image conditioning, using edge detection and pose estimation via canny, hed, and Open Pose Library, and install diffusers, accelerate, and xformers.
Generate images using edges by applying the ControlNet canny edge model with the StableDiffusion pipeline. Tune thresholds and parameters to control edge-based outputs.
Extract edges with the canny edge function from OpenCV and feed only the edges to the neural network to generate high quality images. Experiment with prompts and thresholds.
Extract pose keypoints using the pose detector, then generate photorealistic images with Stable Diffusion via ControlNet, refining with different schedulers and prompts for pose-consistent results.
Generate images from poses by loading internet images, extracting poses with pose_model, and creating pose-driven variations in a diffusion workflow using different schedulers, with yoga examples.
Explore non-xl stable diffusion 1.5 models like epic realism and realistic vision, with practical prompts, negative prompts, schedulers, steps, guidance scales, and correct image dimensions for optimal results.
Explore a google Colab workflow for Stable Diffusion 1.5 and Stable Diffusion XL, including installation, pipelines, prompts and negative prompts, and scheduler selection to generate images.
Explore xl models with stable diffusion, focusing on 1024 by 1024 dimensions and upscaling to boost realism, while using the style selector for diverse predefined aesthetics.
Learn to load and compare stable diffusion models, use prompts and style selectors, adjust image size, and generate high-quality images for logos, art styles, and scenes.
Master ai image generation using stable diffusion covers basics, parameter settings, and initial image generation, plus prompt engineering, custom models, image-guided generation, inpainting, ControlNet, and digital processing techniques.
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The generation of images using Artificial Intelligence is an area that is gaining a lot of attention, both from technology professionals and people from other areas who want to create their own custom images. The tools used for this purpose are based on advanced and modern techniques from machine learning and computer vision, which can contribute to the creation of new compositions with high graphic quality. It is possible to create new images just by sending a textual description: you ask the AI (artificial intelligence) to create an image exactly as you want! For example, you can send the text "a cat reading a book in space" and the AI will create an image according to that description! This technique has been gaining a lot of attention in recent years and it tends to growth in the next few years.
There are several available tools for this purpose and one of the most used is Stable Diffusion developed by StabilityAI. It is Open Source, has great usability, speed, and is capable of generating high quality images. As it is open source, developers have created many extensions that are capable of generating an infinite variety of images in the most different styles.
In this course you will learn everything you need to know to create new images using Stable Diffusion and Python programming language. See below what you will learn in this course that is divided into six parts:
Part 1: Stable Diffusion basics: Intuition on how the technology works and how to create the first images. You will also learn about the main parameters to get different results, as well as how to create images with different styles
Part 2: Prompt Engineering: You will learn how to send the proper texts so the AI understands exactly what you want to generate
Part 3: Training a custom model: How about putting your own photos in the most different environments? In this section you will learn how to use your own images and generate your avatars
Part 4: Image to image: In addition to creating images by sending texts, it is also possible to send images as a starting point for the AI to generate the images
Part 5: Inpainting - exchaning classes: You will learn how to edit images to remove objects or swap them. For example: remove the dog and replace it with a cat
Part 6: ControlNet: In this section you will implement digital image processing techniques (edge and pose detection) to improve the results
All implementations will be done step by step in Google Colab online with GPU, so you don't need a powerful computer to get amazing results in a matter of seconds! More than 50 lessons and more than 6 hours of videos!